# Artistic Stylization

A LoRA trained on the Wan2.1 14B I2V 480p model that transforms any person/object in an image into its Mona Lisa version

LoRA adapter trained on FLUX.1-dev model for text-to-image generation

A LoRA model based on FLUX.1-dev, specializing in text-to-image tasks, particularly suited for generating mandala-style drawings.

This is a text-to-image model based on stable diffusion technology, capable of generating high-quality images from input text descriptions.

A LoRA model for generating anthropomorphic dogs and making random dogs appear in scenes when the trigger word 'HuWoof' is used.

This is a fine-tuned model based on Stable Diffusion XL, specializing in generating silhouette-style images.

This is a Low-Rank Adaptation (LoRA) model based on Stable Diffusion XL, specifically designed to generate highly aesthetic images.

This is an SDXL-based LoRA model capable of transforming any prompt into a watercolor painting effect.

A personalized text-to-image generation model fine-tuned with Dreambooth based on Stable Diffusion v1-5, capable of generating customized image content using the specific concept 'akshyaiD'